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Civitai
civitai.com
Civitai: The Home of Open-Source Generative AI
Explore thousands of high-quality Stable Diffusion & Flux models, share your AI-generated art, and engage with a vibrant community of creators
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Discover Free Stable Diffusion & Flux Models
Browse from thousands of free Stable Diffusion & Flux models, spanning unique anime art styles, immersive 3D renders, stunning photorealism, and more
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Pony Diffusion V6 XL
This model supports a wide array of styles and aesthetics but provides an opinionated default prompt template that allows generation of high quality samples with no negative prompt and otherwise default settings · score_9, score_8_up, score_7_up, score_6_up, score_5_up, score_4_up, just describe what you want, tag1, tag2 · (previous Pony Diffusion ...
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Stable Diffusion Art
stable-diffusion-art.com › home › tutorial › stable diffusion models: a beginner’s guide
Stable Diffusion Models: a beginner's guide - Stable Diffusion Art
September 27, 2025 - Stable Diffusion Models, or checkpoint models, are pre-trained Stable Diffusion weights for generating a particular style of images.
Discussions

Which stable diffusion base models you're using most and for what? (SD 1.5, SDXL, or other stable diffusion models)
pony for nsfw, flux for sfw. i have ~10 pony models i run a prompt through to pick what i like best, mix of realism + anime ones. sd1.5/sdxl are also great, tons of loras and really artistic checkpoints, just kinda moved on from them because of their inconsistency with hands/feet and difficulty doing good nsfw. More on reddit.com
🌐 r/StableDiffusion
4
2
September 14, 2024
What’s the best Stable Diffusion models out there?
Very broad question. Very very very broad question. For example, if you want to make depth images without a reference image, then depthstyle is the best model. If you want general recs, i'll give you a list of my most used models. RMSDXL Orion, Aries, Corvus, Scorpius. These are my versatile workhorse models. I generally always start with these, moving on to the others if i want something specific. AnimaginexlV3, best XL anime model. Good prompt adherence. Cheyenne_v14, just started using it but it's awesome western comicbook style. Envy Gacha, if you like that over the top fashion anime style. JuggernautXLv8, i like this one for close ups. realstockphoto_v2 combined with bad quality lora for realism. Those are all XL, I occasionally use 1.5 but i'm not as well versed with them. I mostly use 1.5 for LORAs. Mistoon anime has that screenshot of a tv show look. 3danimationdiffusion for a pixar vibe. Juggernautaftermath, wafflemix, realismfromhades for realism. And then several more you'd have to find on your own. Note some of these might not be the "best" but they do the job i want them to do. Took a lot of hoarding and testing to figure that out though, so good luck with it. More on reddit.com
🌐 r/StableDiffusion
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31
January 30, 2024
Which Stable Diffusion should use? XL, 3.5 or 3.0?
So here's the history SD 1.4 -- the first model launched in August 2022 . . . obsolete now, no one uses it SD 1.5 -- the first widely adopted model. Still useful today, particularly useful on less powerful computers, low memory GPUx; models were typically 2 gigabytes or so. Training and render size was small, typically 640 x512, typically, so you need to use upscalers. Still useful in particular cases, its fast, its small and in particular use cases its actually best. ControlNets seem to work better with SD 1.5 than SDXL or Flux . . . SD 2 & 2.1 -- basically failures. Not much used, no reason to SDXL -- a big jump up, trained on and generates natively in 1024 x 1024. Models are typically 6 gigabytes (FP 32 versions of these do exist, useful in training new checkpoints, but not for inference). I still use SDXL a lot, usually on the Fooocus UI. There are also two SDXL derivatives worth mentioning, Pony and now Illustrious. These are based on SDXL, but have been tweaked enough that they respond to prompts quite differently, have their own LORAs and so on. The big advantages of these models is that they respond to text prompting better than base SDXL does, particularly for posing . . . if you're trying to describe, say, acrobats or athletics, Pony or Illustrious models will likely respond better to something like "orc grabbing kangaroo by the tail and swinging it". Pony and Illustrious both started out as anime/manga oriented models, but there are realistic variations. Also in the SDXL family, there's SDXL Turbo and Lightning, designed for speed, which make posslble realtime implementations (eg you draw something on the screen and you get an AI'd version of the brush stroke -- Kreia is a great implementation of this, and you can get some implementations in ComfyUI on the desktop, if you have a good GPU). SD 3 -- failure, little used. SD 3.5, much better than 3 fixed some of its glaring errors. If it were competing with SDXL, then there might be some interest . . . but its not . . . the state of the art is FLUX -- built by Black Forest Labs, the team that originally developed Stable Diffusion, FLUX is state of the art. In the "Schnell" variant it can be very quick, in the "Dev" variant, at the cost of slow speed and huge models (11 gigabytes and up), the quality is excellent, rivaling and exceeding Midjourney. Training LORAs is incredibly easy with FLUX Gym . . . for most users, IFF they have an adequate GPU (Nvidia 3000 or 4000 series GPU, preferably with lots of VRAM) . . . this is the shortest path to "looking like a photograph" . . . with that said, a skilled user who's willing to dig into tools and techniques can produce beautiful images from that SD 1.5 technology, now 2.5 yrs old. More on reddit.com
🌐 r/StableDiffusion
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March 24, 2025
So what is the best model to use locally now? And how to do it the simplest way?
I've been mainly focusing on Flux for the last few months. It's great at realistic stuff, but struggles with styling and needs powerful hardware. Forge is probably the easiest option with good support for Flux but I think it might be falling behind on things like ControlNet. I still think ComfyUI is worth the effort to learn since it's the most powerful option and supports pretty much everything. More on reddit.com
🌐 r/StableDiffusion
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December 26, 2024
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Reddit
reddit.com › r/stablediffusion › what’s the best stable diffusion models out there?
What’s the best Stable Diffusion models out there? : r/StableDiffusion
January 30, 2024 - I would suggest visiting civitai.com they rank models according to most downloads, popularity, etc. And there are lots of reviews which can help. Remember that there are different types of Stable Diffusion Models, namely Base Models, like SDXL, Pony and there are also trained models and merged models.
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Hugging Face
huggingface.co › models
Models – Hugging Face
Docker Model Runner · Lemonade · Inference Providers Select all · Groq · Novita · Nebius AI · Cerebras · SambaNova · Nscale · fal · Hyperbolic · Together AI · Fireworks · Featherless AI · Zai · Replicate · Cohere · Scaleway · Public AI · OVHcloud AI Endpoints · HF Inference API · WaveSpeed · Misc Reset Misc · stable-diffusion ·
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Wikipedia
en.wikipedia.org › wiki › Stable_Diffusion
Stable Diffusion - Wikipedia
2 days ago - Stable Diffusion is a latent diffusion model, a kind of deep generative artificial neural network. Its code and model weights have been released publicly, and an optimized version can run on most consumer hardware equipped with a modest GPU ...
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GitHub
github.com › CompVis › stable-diffusion
GitHub - CompVis/stable-diffusion: A latent text-to-image diffusion model
Stable Diffusion is a latent text-to-image diffusion model. Thanks to a generous compute donation from Stability AI and support from LAION, we were able to train a Latent Diffusion Model on 512x512 images from a subset of the LAION-5B database.
Starred by 72K users
Forked by 10.5K users
Languages   Jupyter Notebook 88.8% | Python 11.1% | Shell 0.1%
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AWS
aws.amazon.com › what is cloud computing? › cloud computing concepts hub › generative ai › what is stable diffusion?
What is Stable Diffusion? - Stable Diffusion AI Explained - AWS
2 days ago - Stable Diffusion is a generative artificial intelligence (generative AI) model that produces unique photorealistic images from text and image prompts. It originally launched in 2022. Besides images, you can also use the model to create videos and animations. The model is based on diffusion ...
Find elsewhere
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Stability AI
stability.ai › stable-image
Stability AI Image Models — Stability AI
Advanced text-to-image generative model with 3.5 billion parameters, offering high-resolution and photorealistic outputs. ... A distilled version of SDXL, utilizing Adversarial Diffusion Distillation to enable real-time image generation in as ...
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Stable Diffusion
stablediffusionweb.com
Stable Diffusion Online
Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, cultivates autonomous freedom to produce incredible imagery, empowers billions of people to create stunning art within seconds.
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YouTube
youtube.com › watch
Stable-Diffusion: How To Pick Your MODEL (In 2 Minutes!) - YouTube
Here's how you find and apply a new base model when training your own images -CivitAI Link Here:https://civitai.com/modelsIf you enjoyed this video, please c...
Published   March 7, 2024
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Ikomia
ikomia.ai › home › blog › top ai diffusion models: ultimate comparison & guide [2024]
Top AI Diffusion Models: Ultimate Comparison & Guide [2024]
September 10, 2024 - This makes it exceptionally useful for creating photorealistic portraits or fashion illustrations that require a unique, distinct finish. ... RealVisXL V4.0 stands as the top Stable Diffusion model for crafting lifelike human images.
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Medium
medium.com › @onkarmishra › stable-diffusion-explained-1f101284484d
Stable Diffusion Explained. How does Stable diffusion work… | by Onkar Mishra | Medium
August 21, 2024 - In latent diffusion, the model is trained to generate latent (compressed) representations of the images. ... Stable Diffusion is a large text to image diffusion model trained on billions of images. Image diffusion model learn to denoise images ...
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Openxcell
openxcell.com › home › 15 best stable diffusion models to create lifelike ai visuals
Best Stable Diffusion Models for Hyperrealistic AI Art
April 14, 2025 - Deliberate is a versatile and balanced model which merges realism with artistic flair. It makes clean, expressive images that retain a natural yet stylish look. ... Dreamshaper is a fan favorite because it blends realism with artistic imagination, making it perfect for both realistic and creative outputs. If you are wondering what is the best stable diffusion model for flexibility across styles, DreamShaper is a top choice.
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Secoda
secoda.co › glossary › what-are-stable-diffusion-models
What Are Stable Diffusion Models? - Explanation & Examples | Secoda
A Stable Diffusion model is a type of artificial intelligence (AI) model designed for generating digital images from textual descriptions.
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Hugging Face
huggingface.co › CompVis › stable-diffusion
CompVis/stable-diffusion · Hugging Face
Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. This model card gives an overview of all available model checkpoints.
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Aiarty
aiarty.com › stable-diffusion-guide › best-stable-diffusion-models.htm
40+ Best Stable Diffusion Models 2025 (Free & Updated)
1 week ago - Stable Diffusion has become the go-to text-to-image AI tool, letting creators generate stunning visuals from text prompts. With thousands of Stable Diffusion models available, each tailored to a specific style, it can be overwhelming to pick the right one. This guide – updated for 2025 – highlights the best Stable Diffusion models, new hardware-optimized releases like SD 3.5 Large, and next-gen models like Flux 1.1 Pro.
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Medium
medium.com › @promptingpixels › comparing-stable-diffusion-models-2c1dc9919ab7
Comparing Stable Diffusion Models | by Prompting Pixels | Medium
April 2, 2024 - Comparing Stable Diffusion Models Stable Diffusion, an open-source text-to-image model released by Stability AI, has revolutionized the field of generative AI. Since its initial release in 2022 …
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GitHub
github.com › civitai › civitai › wiki › How-to-use-models
How to use models
How you use the various types of assets available on the site depends on the tool that you're using to generate your images using Stable Diffusion. We've included instructions for a few of them here, but if your tool isn't here, we invite you to open a Q&A Discussion Post providing directions so we can include it here. ... Make awesome images! ... Make awesome images! ... Place the aesthetic graident inside the aesthetic_embeddings directory of your AUTOMATIC1111 Web UI instance ... Make awesome images! ... Place the hypernetwork inside the models/hypernetworks directory of your AUTOMATIC1111 Web UI instance
Author   civitai
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Ommer-lab
ommer-lab.com › home › research › high-resolution image synthesis with latent diffusion models
High-Resolution Image Synthesis with Latent Diffusion Models - Computer Vision & Learning Group
July 12, 2023 - By decomposing the image formation process into a sequential application of denoising autoencoders, diffusion models (DMs) achieve state-of-the-art synthesis results on image data and beyond. Stable diffusion.
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Stable Diffusion XL
stablediffusionxl.com › home › stable diffusion › stable diffusion xl – download sdxl 1.0 model here
Stable Diffusion XL - SDXL 1.0 Model - Stable Diffusion XL
July 17, 2024 - Stable Diffusion XL (SDXL), is the latest AI image generation model that can generate realistic faces, legible text within the images, and better image composition, all while using shorter and simpler prompts.